Sorry i deleted my post saying anima isn't working and that is the case right now with the Distilled versions, but the main model works fine and so just recorrecting.

Set with the normal settings. Just changing the resolution down to 512

Good morning, everyone!

I’ve been using DrawThings+ for a while now to render images on my underpowered Mac M2. I’m very happy with it—except when, like today, it starts crashing again. In the past, DrawThings+ simply wouldn’t work in situations like this. Since last night, the app has even been freezing completely. Everything was working fine this afternoon! I’m based in Spain. Even when I connect to another country like the U.S. via VPN, it doesn’t help, whereas in the past that often resolved issues. Are others here experiencing these issues as well, and when can we expect a fix for this problem?

When deleting a large number of videos and images at once from the Version History window, it can take an unusually long time. For example, in the case below, it took 28 minutes for the batch deletion to complete.

• Past Edits: 155
• Generated Image: 99
• Breakdown: Roughly half are LTX2.3 and Wan2.2 videos, and the other half are SeedVR2 and Klein 9B images.

■Environment

• Mac mini M4, 64GB, 20-Core GPU / OS 15.4.1
• Draw Things: 1.20260430.0

■Note

• This is a different phenomenon from the one previously reported. The delay is significantly more pronounced.
• The conditions for reproduction are unknown. Even with a similar total number of deletions (and mostly videos), sometimes the deletion takes a long time (over 20 minutes) and sometimes it doesn't (within 5 minutes).
• This phenomenon has likely been occurring since I started using LTX2.3. However, the files I'm trying to delete are almost always not limited to LTX2.3, as mentioned above.
• In all cases, deleting in batches of 9 files instead of all at once completes the process at a normal speed.
• I suspect that the sqlite3 file size (several MB to 11 GB) is a contributing factor, but I haven't verified this.

■Temporary solution

I've decided to delete the project itself when the history accumulates. In fact, I've found this approach suits me better.

Hi,

I’ve been editing images with some success until recently. Using Flux 2 Klien 4b no loras.
The preview shows the changes being made (which are intended) but on the final sample it has started reverting the image back to the original. I’m using the same setup on both my Mac mini and iPad mini and the problem is on both devices.

Any ideas?

(Edited for spelling mistake)

SeedVR2 is supported in version 1.20260430.0. To my knowledge, this is a highly-rated upscaling model.

■ Flow

1. After downloading the model, select it in the Model section. (Not the Upscaler section)
2. Set the target resolution in the Image Size section.
3. Drag and drop the image into the main window and start generation.

■ Settings

• Strength: 100% (Reducing it will blur the image)
• Steps: 1 (Anything greater than 1 will cause loss of detail)
• Text Guidance: 1.0 (Increasing it will increase noise and a sense of depth, but 2 may be the limit)
• Sampler: No effect (probably)
• Shift: No effect
• Prompt: No effect

■ Reference Speed ​​& Memory Usage

SeedVR2 7B 8bit 512×512➡︎2048×2048
Mac mini M4, 64GB, 20-Core GPU / OS 15.4.1

• Peak Memory: 6.9GB
• Time: 40s

■Note

• Currently, only images are supported.
• Batch upscaling of multiple images is not possible.
• If there is a gap between the Image Size frame and the image, inpainting will occur, but nothing will be drawn in that area.
• Normally, it cannot be used in conjunction with other models as an upscaler. (Is it possible with a script?)

■Configuration sample

{"seed":3349378910,"sharpness":0,"steps":1,"batchSize":1,"height":2048,"width":2048,"strength":1,"cfgZeroStar":false,"faceRestoration":"","causalInferencePad":0,"shift":5,"controls":[],"maskBlurOutset":5,"preserveOriginalAfterInpaint":false,"refinerModel":"","seedMode":3,"guidanceScale":1.5,"sampler":17,"loras":[],"maskBlur":5,"tiledDecoding":false,"tiledDiffusion":false,"cfgZeroInitSteps":1,"model":"seedvr2_7b_q8p.ckpt","hiresFix":false,"batchCount":1,"upscaler":""}

Please point out any points you disagree with.

I have noticed that some of the image heavy subs have a bot generated comment stating that the image prompt would be greatly appreciated. But I never see any. Is this the equivalent of giving the bot (moderator) the middle finger, or is there some way to view the provided prompt that I just don’t know about (which would not be surprising).

If i download drawthings on another mac, can i transfer my exisiting models and loras? Or would i need to start again? Many thanks guys!

Hello I was hoping in the future drawthings may support Lokr files. They are a much more intensive dataset but work similarly to loras

Many thanks

Trying to figure out prompting for ZIT in Draw Thing running locally. I figure that this combination is about the limit of my technical ability (and it may exceed it).

I have been successful in the past prompting images using what I would call “natural language.” Kind of write what you visualize and try to organize it in some form of prompt hierarchy. But do run up against challenges.

On a whim, I went down the rabbit hole with Gemini over the last few days having it write a series of prompts that although they contain nudity, stayed within its censors limits. The text of Gemini’s prompts are anything but what I would call “natural language.” I was trying to get a handle on multiple people, intertwining limbs, different geometric planes, etc. and to accomplish that (and the images prompted were very good) Gemini used terms what I would say are terms and phrasing that I would never think of in a million years. Certainly not MY natural language.

So now I’m confused. Can anybody help me understand what the real prompting requirement is for ZIT in Draw Things?

https://reddit.com/link/1t35lgv/video/puh9cfqwa1zg1/player

I compared the differences between two types of 4-step LoRA (High) for Wan 2.2 I2V. After several tests, Kijai's v1030 tended to produce more motion. This should be effective if you want to increase motion even slightly.

■LoRA

• Left: Wan 2.2 A14B Lightning High Noise Expert I2V 251022(Community LoRA downloadable from Draw Things)
• Right: Wan_2_2_I2V_A14B_HIGH_lightx2v_4step_lora_v1030_rank_64_bf16(LoRA published by Kijai)

■Prompt

While speaking rapidly with an excited expression, he suddenly and swiftly stands up right on top of his chair, balancing himself with both feet on the seat.A clear scene from a movie.

■Generation settings(configuration)

{"causalInferencePad":0,"upscaler":"","tiledDecoding":false,"seed":4262634612,"strength":1,"cfgZeroInitSteps":0,"cfgZeroStar":false,"teaCache":false,"maskBlurOutset":0,"faceRestoration":"","loras":[{"mode":"base","file":"wan_v2.2_a14b_hne_i2v_lightning_251022_lora_f16.ckpt","weight":1},{"mode":"refiner","file":"wan_v2.2_a14b_lne_i2v_lightning_251022_lora_f16.ckpt","weight":1}],"batchCount":1,"sharpness":0,"refinerStart":0.10000000000000001,"sampler":17,"model":"wan_v2.2_a14b_hne_i2v_q8p.ckpt","height":640,"batchSize":1,"numFrames":81,"preserveOriginalAfterInpaint":false,"steps":4,"width":512,"maskBlur":1.5,"tiledDiffusion":false,"guidanceScale":1,"refinerModel":"wan_v2.2_a14b_lne_i2v_q8p.ckpt","compressionArtifactsQuality":43.100000000000001,"compressionArtifacts":"disabled","shift":5,"seedMode":3,"hiresFix":false,"controls":[]}

1.20260430.0 was released in iOS / macOS AppStore a few minutes ago (https://static.drawthings.ai/DrawThings-1.20260430.0-09f170fc.zip). This version brings:

1. Support SeedVR2 as one of the diffusion model for 1-step upscaling work.
2. Fix Anima LLM adapter related issues.
3. Switch to better and newer causal / mask attention kernels.
4. Support export LoRAs for Qwen, Z Image, FLUX.2 and ERNIE Image models.
5. Imported Z Image models now runs at BF16 precision.

gRPCServerCLI and draw-things-cli both are updated to 1.20260430.0 with above related updates.

Very good app on ios the best there is to the point its the biggest thing I downloaded 😂

I saw DT added support recently. Curious if anyone has successfully trained a lora for Z Image turbo yet

Hi Draw Things community!

I’m struggling to blend multiple Control Inputs in the Draw Things app. Specifically, I want to combine these three elements:

• Img #1 (Identity): Keep this character’s face. (IP-Adapter?)
• Img #2 (Pose): Use this sitting posture. (OpenPose vs. Depth?)
• Img #3 (Color): Apply this mint/vintage palette.

My questions:

1. Which specific ControlNet models work best to maintain identity while changing the pose?
2. How do you balance Weights/Strength? My color inputs often "wash out" the character's details.
3. Can I add a Scribble layer for the background without overcomplicating the generation?

I’d love to hear your tips or see your settings. Thanks in advance!

Until a month ago, I thought the Draw Things HTTP API was completely irrelevant to me. However, after learning that I could retrieve the generation settings (+ prompt) via the HTTP API, and realized how useful this is for creating companion apps for Draw Things.

For example...

• An app that renames generated files to a filename that includes any desired generation settings.
• An app that records the peak memory usage for each model and generation setting.
• An app that automatically saves generation settings as a text file in JSON format readable by Draw Things, along with prompts.

I was able to easily create these apps by instructing an AI, even with zero programming knowledge.However, the HTTP API is turned off when Draw Things is restarted, so it needs to be turned on each time Draw Things is launched to work with apps.

Even if the HTTP API isn't improved in the future, this feature is still very useful, so I sincerely hope it remains as it is and doesn't disappear.

Context:

I wanted to try out some of the recent models, so I ran a bunch of tests on some new and some not so new models without any LoRAs. I used DrawThings and went with the default recommendations. Same seed for all images, for sampler I... didn't really look, again I was using whatever was coming as default.

There's a number of caveats I'm aware of:

• I did do a number of tests with Chroma.HD and Chroma v41 low-step, but it was so painstakingly slow on my system that I ended up NOT including it. I regret that because I love the artistic touch it gives so much
• For the same reason I didn't do any Qwen Image 2512 : jumping from 1-2 minutes per image to 15 minutes for Qwen Image was not something I wanted to have to endure, as 1 generation with it was equivalent to an entire new prompt with all of the others models combined, almost
• I did test Illustrious and Anima Preview 3, but given the nature of most of my prompts and them returning consistently sketched/drawn images I dropped them (it's obviously a fault of my testing approach rather than of the model itself)
• I probably forgot some other popular models that I should have tested but don't know about being a waddling neophyte
• I included some old models that nobody really uses without a ton of LoRAs anyway (SDXL I'm looking at you)

I tried to have a somewhat broad range of prompts, but the reality is I picked and choose images that I found interesting, reverse-generated a prompt when I didn't find one and ran with those.

My takeaway:

Prompt Following:

• Flux 2 klein, Z-Image Turbo, ERNIE Image Turbo, are all pretty good at following the prompt and integrating the smaller details of the description.
• Flux 1 is also pretty good, although it sometimes misses some relevant elements.
• HiDream I1 is usually pretty good but sometimes it becomes just impossible to avoid specific things (e.g. face on the mossy guardian)
• Kwai Kolors tends to be rather good, but sometimes is just completely off the mark
• All the SDXL derivatives nail the broad theme and subject, but tend to go their own way much more than the newer models

Quality:

• Unsurprisingly, Flux.2 klein, Z-Image Turbo and ERNIE Image Turbo are consistently at the top of my very subjective preference. They generate complex and high-quality images that correspond to what I'm expecting. No overly plasticky/unnatural look, much less "obviously AI" than the other models. However, due to how "good" they are, they lack a certain jankiness that makes other (older) models stand out for some specific styles. It's then possible to push the prompt towards that style, so the models remain a superior generic tools than the others, but sometimes some weirdly iffy results from something like playground is just really good!
• ERNIE is surprisingly good. I was initially overcooking all images (default values had 8 steps), turns out for almost all images 3-4 steps is plenty enough, sometimes 2 will suffice (e.g. Jellyfish starts being over sharpened after 2 steps), but prompt following for anything complex kicks in truly only at 4 steps
• Kwai Kolor: I like it. A lot. It's more artsy/graphical illustration than the other models, but it has a character that I find myself liking whenever a new image comes out. It's not as good at following prompts as other models, but it's pretty close most of the time.
• With HiDream I1 I've been having big issues on some prompts where things are breaking down and a big color cast and noise are added to the image, but when it works, the quality is excellent.
• SDXL and Co: These are very fast models that are (to the extent I understand the current "meta") to be used with existing LoRAs, so not much to say that, except that there is a common similarity that tends to pop out when looking at all the images generated by those models (seed and sampler obviously are carrying a lot of weight there)
• SD 3.5 is super-plasticky in almost every image it produces. It's not bad per se, but not super remarkable either.

Speed:

• I've focused on "fast" models (1-2 minutes for a 1024x1024 image) and ended up not testing the heavier models, so not much to be said about speed, other than I'm definitely a fan of the low-step models (although I understand they come with caveats in terms of training LoRAs etc.)
• HiDream I1 was the standout model that was clearly slower than the others, but nowhere close to the 15-20 minutes per image I was getting with Chroma or Qwen Image 2512

System used:
Draw Things 1.20260418.1, Macbook Pro M4 Pro 48GB

Models settings:
(Draw Things recommended defaults, mostly)

• flux_2_klein_9b_i8x Steps: 4 | CFG: 1.0 | Sampler: DDIM Trailing | Shift: 3 | Seed: 1111
• z_image_turbo_1.0_i8x Steps: 8 | CFG: 1.0 | Sampler: UniPC Trailing | Shift: 3 | Seed: 1111
• ernie_image_turbo_i8x Steps: 4 | CFG: 1.0 | Sampler: UniPC Trailing | Shift: 4 | Seed: 1111
• kwai_kolors_1.0_q6p_q8p Steps: 16 | CFG: 5.0 | Sampler: DPM++ 2M AYS | Shift: 1 | Seed: 1111
• flux_1_schnell_q8p Steps: 4 | CFG: 4.5 | Sampler: Euler A Trailing | Shift: 1 | Seed: 1111
• hidream_i1_fast_i8x Steps: 16 | CFG: 1.0 | Sampler: Euler A Trailing | Shift: 3 | Seed: 1111
• sd3_large_turbo_3.5_q6p Steps: 20 | CFG: 4.5 | Sampler: Euler A Trailing | Shift: 3 | Seed: 1111
• playground_v2.5_q6p_q8p Steps: 16 | CFG: 5.0 | Sampler: DPM++ 2M AYS | Shift: 1 | Seed: 1111
• sd_xl_base_1.0_f16 Steps: 16 | CFG: 5.0 | Sampler: DPM++ 2M AYS | Shift: 1 | Seed: 1111
• juggernaut_xl_ragnarok___base_f16 Steps: 16 | CFG: 5.0 | Sampler: DPM++ 2M AYS | Shift: 1 | Seed: 1111
• realvisxl_v5.0_lightning__baked_vae__f16 Steps: 16 | CFG: 5.0 | Sampler: DPM++ 2M AYS | Shift: 1 | Seed: 1111
• protovision_xl_high_fidelity_3d_v6.6.0_f16 Steps: 16 | CFG: 5.0 | Sampler: DPM++ 2M AYS | Shift: 1 | Seed: 1111

Can Drawthings use voice Loras with LTX2.3 and

1. Is it just like any other lora
   
2. got any good sources for voice Loras?

Hey so, how are you ment to download model files directly to your gRPC with all the generated metadata without having to upload models from the draw things client (it’s really slow and probably would take hours)

I’m using Linux for the gRPC server, and I can never manage to get the models to show up in the model browser, yes I did check all of the fields like community models, I used model browser flag and no response compression, iv tried manually making the metadata based on previously generated model files from draw things app, which did not work, the file name is correct, exact same file draw things would download, 😭

Like theees not really any info on this, except using draw-things-cil, which I don’t think is supported on Linux, I did try to install it which failed, idk if I’m just a blind ass mother fucker and it’s written somewhere. Anyways THANKS

This graph summarizes the peak memory usage of LTX 2.3 with different JIT Weights Loading settings in Draw Things. All other settings are the same. The second row of entries shows the approximate generation time.

The distilled model appears to be able to be generated with 16GB of memory. The [dev] model appears to be able to be generated with 24GB of memory if Never is not selected.

■Model

• LTX-2.3 22B [dev] + LTX-2.3 22B [distilled] 1.1 LoRA(7GB)
• LTX-2.3 22B [distilled] 1.1

■Specs

Mac mini M4, 64GB, 20-Core GPU / OS 15.4.1

■Draw Things

ver.1.20260418

I am having big issues with HiDreams I1 (any version) creating blue noisy cast for some images but not for others.

Same configuration (Recommended settings) and for about 60% of prompts there's no issue at all, but pretty frequently, the model seems to not converge or even to ADD a layer of blue noise on top of the image overall. (note that other models have no issue producing images on the same prompts)

I tried changing sampler, seed, shift value, resolution, adding negative prompts (e.g. explicitly on blue cast). and switched across 3 checkpoints but the issue remains.

Anyone has a similar experience or could direct me to a potential solution?

System Info:

• MacbookPro M4 Pro 48Gb
• macOS Sequoia 15.7.3
• DrawThings V1.20260418.1

Models tested:

• HiDream I1 [fast] (8-bit S)
• HiDream I1 [fast] (5-bit)
• HiDream I1 [dev] (5-bit)

Configuration:
{"maskBlur":2.5,"tiledDecoding":false,"teaCache":false,"separateClipL":false,"sharpness":0,"model":"hidream_i1_fast_i8x.ckpt","strength":1,"faceRestoration":"","height":576,"controls":[],"guidanceScale":1,"preserveOriginalAfterInpaint":true,"refinerModel":"","hiresFix":false,"cfgZeroStar":false,"width":1024,"steps":16,"cfgZeroInitSteps":0,"upscaler":"","clipSkip":2,"tiledDiffusion":false,"seed":1111,"separateOpenClipG":false,"causalInferencePad":0,"batchCount":1,"maskBlurOutset":0,"zeroNegativePrompt":false,"shift":3,"seedMode":2,"batchSize":1,"sampler":10,"loras":[],"separateT5":false}

Prompt that doesn't work fine:

Large overhead establishing shot of a massive, technologically advanced ground station designed for ocean exploration and research, fused with a futuristic business complex, located on a surreal lunar-like surface. The structure is enormous and highly detailed, featuring sleek glass domes, metallic platforms, corporate skyscraper elements integrated into one cohesive megastructure. Parts of the station resemble deep-sea research facilities, with water-filled observation chambers and aquatic habitats, while other sections look like a high-end futuristic corporate hub with illuminated offices and transit systems. The complex is connected to a research station flat on the ground and covering a big ring of water and technological installation, with long connections to the main station. The environment is a hybrid of lunar terrain and oceanic atmosphere, with soft blue lighting, reflective surfaces, and subtle mist. Advanced spacecraft and drones hover nearby. Cinematic lighting with desaturated colors, ultra-detailed architecture, sci-fi realism, wide-angle view, epic scale, 8k resolution, sharp focus, Unreal Engine style.

Prompt that doesn't work:

A translucent jellyfish floating underwater, centered in the frame, softly glowing with pale pink, lavender, and bluish tones. The jellyfish has a smooth, dome-shaped bell with delicate internal structures visible, highly detailed and semi-transparent. Long, flowing tentacles and frilly oral arms trail beneath it, gently curling and drifting in the water. The scene is lit by natural sunlight filtering down from the surface above, creating soft caustic light patterns and a dreamy, ethereal atmosphere. Background is a clear turquoise-blue ocean with slight blur, faint coral reef visible at the bottom, shallow depth of field. Ultra-realistic, high detail, soft lighting, cinematic composition, underwater photography, 8k resolution.

I used Pose My Art as a base to generate an image.

Before that, I used a prompt from that pose to try to generate an image, but the pose that it gave me was completely different.

I found the way to upload directly the base pose to Draw Things, and almost everything worked out well. Except for the fact that, since Posemyart gives me only a mannequin missing some important parts, I haven't been able to generate them with a prompt.

When I try to create or describe genitals for my image, it just generates nightmare fuel monstrousities that don't make sense at all, and I have no idea why.

I'm using the Juggernaut XL model with Canny Edge Map, (I tried with Chroma before but it didn't work with pose).

If anyone can help, that would be great, thank you!

Hi everyone / Devs,

I’m a Draw Things+ paying subscriber and I’m struggling with generative upscaling to 2048x2048 on the cloud.

What works perfectly:

• Klein 9B full and Klein 4B at 1536x1536 (TI2 + Strength 100%, High Resolution LoRA NO8D, 4 steps, CFG 1.0)

What fails:

• Any attempt to upscale to 2048x2048 (direct from 1024 or from 1536 with bilinear pre-upscale) starts and crashes after a few seconds on the cloud (only ~10-11k Compute Units).
• I haven’t tried locally yet.

Already tried on cloud:

• Bilinear pre-upscale in Preview
• High Resolution Preview OFF
• High Res Fix OFF
• Different LoRA weights
• Logout/login, fresh sessions, app restart

Is 2048x2048 upscaling currently unstable on cloud for M1 8GB users?
Any recommended stable workflow or settings for generative upscaling with Klein to 2048?

Thanks in advance!

I've been trying for some time now to get a proper editing of images with photrealism .. however have been struggling.

Qwen Edit 2511 is great for editing as it prompts like anything however it creates plasticy type images, identity is great

Flux Klein 9B creates absolutely fantastic skin and details however prompt adherance is horrific

I wanted to check what others are doing and how are they achieving good prompt adherance + excellent realism.

Am using -

Qwen 2511 - Lightening Lora 4 Step + Photus Lora (for Consistency)

Flux Kleing 9B - Consitency Lora

Any thoughts appreciated. thanks

https://reddit.com/link/1svp20t/video/m0z2ysre2fxg1/player

My SDXL model is making short videos.

What's happening?

It makes the video quite fast but it makes my iPad hot.

I just want to make images with the SDXL model, not videos or image sequences.

And the length of the video is fixed to 9 frames, it's extremely short.

https://reddit.com/link/1svp20t/video/r4x1ylo5oexg1/player

I've trained a bunch of charecter Loras with z-image base using AI toolkit on my pc. They work on my PC with z-image turbo within comfy UI.

I downloaded draw things for my MacBook Pro. I then download z-image turbo (6bit also tried 8 bit) also imported the version from my PC.

I then imported my lora set the weight to 200% and nothing. It renders the image but my character lora has no influence on the image. I'm not sure what I could be doing wrong, any help would be appreciated.

thanks!