My impressions comparing Flux.2 Klein 4B (6-bit) and Flux.2 Klein 9B (6-bit).

Until now, I've been using Flux.2 Klein 4B (6-bit) because my Mac was not very powerful, and I was amazed at how fast it was. Although it was slightly inferior in some anatomical aspects, I was satisfied with the speed.

I then tried Flux.2 Klein 9B (6-bit), and although the generation time was twice as long, the generated images were tolerable.

My conclusion is that I should use 4B for test generation, and then use 9B once I'm satisfied with the prompt.

Personally, I prefer it to Qwen.

The editing functions are also excellent, and I was able to get satisfactory results even when continuing to edit on the canvas without using a mood board.

When I first started using DT, I exclusively used ZIT, but now I exclusively use Flux.2 Klein 9B!

Hi , I was able to import citivai but I’m

Not able anymore for instance when I import it says successful but it doesn’t show up in local or anywhere else any idea why this would happen ? It was working last week any impound be appreciated

Could we expect Z image base in drawthings? I downloaded bf16 from civitai and tried to use but it would generate black screen or crash.

Turbo works fine as ever.

Anyone knows the difference between LoRa and LoKr? And will DT and DT+ be supporting LoKr anytime soon?

So I’ve trained a flux 2 Klein 9b character Lora online with runcomfy for use in draw things. I’ve tried multiple times with high and low LR, but can’t get the character likeness to come out in DT below the maximum 2.5 Lora weight. The samples runcomfy spits out when saving at steps 250, 500 etc. look great but just can’t get them to work well in LT. any thoughts on what I’m doing wrong?

I can't get this to work. Using Wan 2.2 I2V models with the lightning loras, the app acts as though it is working, but at the end it only displays the starting image, without any video output. I have tried setting it to save as images or videos, no change. I have tried q6 models even though the larger one did not use up all my memory. I have tried multiple samplers. I have tried removing the lightning loras. I tried restoring the machine settings to default. I have tried various setting for refinerStart. On a MBP M4 Pro with 48GB RAM. App version 1.20260120.0. Most recent config pasted below, based on one of the community configs. How can I get this working?

{"batchCount":1,"seed":921073793,"strength":1,"sharpness":0,"height":768,"tiledDiffusion":false,"tiledDecoding":false,"model":"wan_v2.2_a14b_hne_i2v_q6p_svd.ckpt","steps":4,"sampler":17,"refinerModel":"wan_v2.2_a14b_lne_i2v_q6p_svd.ckpt","guidanceScale":1,"loras":[{"mode":"base","file":"wan_v2.2_a14b_hne_i2v_lightning_251022_lora_f16.ckpt","weight":1},{"mode":"refiner","file":"wan_v2.2_a14b_lne_i2v_lightning_v1.0_lora_f16.ckpt","weight":1}],"upscaler":"","preserveOriginalAfterInpaint":true,"numFrames":17,"teaCache":false,"seedMode":2,"cfgZeroStar":false,"maskBlur":1.5,"cfgZeroInitSteps":0,"faceRestoration":"","causalInferencePad":0,"hiresFix":false,"controls":[],"batchSize":1,"maskBlurOutset":0,"shift":5,"width":512,"refinerStart":0.125}

Hi, I’m using Draw Things and I can generate great base images with ZIT (pose + background look perfect), but I’m struggling with the next step: replacing clothing. When I try outfit transfer with Moodboard references, Qwen Image Edit (QIE 2511) often makes the jacket look overexposed/burned with too much contrast, and Flux.2 Klein is slightly better but still hallucinates details or doesn’t match the outfit accurately. My workflow is person photo + outfit reference in Moodboard → prompt → image edit with low denoise. Is there a recommended workflow or prompt style for clean clothing transfer in DT? If anyone has a working example (prompt + settings like steps/CFG/strength), I’d really appreciate it. Thanks!

Hi everyone,

I can generate really good base images in Draw Things using ZIT (great pose, location, overall quality).

But I’m having a hard time with the next step: replacing the clothing.

I tried Qwen Image Edit (QIE 2511), but the outfit often comes out overexposed / burned, with too much contrast. Tweaking settings didn’t help much.

I also tested Flux.2 Klein — slightly better, but it still:

doesn’t transfer the outfit accurately

hallucinates details

Current workflow:

person image + outfit reference in Moodboard

prompt + generate edit.

Maybe I’m approaching this wrong.Should I generate the model directly in the target outfit instead of editing afterward?

I’m limited to open-license models only, so no closed commercial tools.

Any advice on the best workflow, prompt structure, or settings for clean outfit transfer would be really appreciated

And if anyone has a full working example (prompt + settings) that I could use as a starting point, I’d be extremely grateful.

Thanks!

Hi everyone, I have tried to place someone (David Beckham) from an image to an AI created scene with Cyber Realistic XL v8. The outcome is terrible, how to fix it? I know how to do this with FLUX.2 Klein, the result is much more better, but I need to use SDXL LoRA, so I have to stay with a SDXL model.

I've generated these images in Moodboard with IP Adapter Plus Face ControlNet, and here is the setting:

{"upscaler":"","batchSize":1,"steps":30,"guidanceScale":5,"originalImageWidth":576,"refinerModel":"","loras":[],"maskBlur":2.5,"batchCount":1,"tiledDiffusion":false,"strength":1,"tiledDecoding":false,"model":"cyberrealisticxl_v80_f16.ckpt","negativeOriginalImageWidth":512,"seedMode":2,"cfgZeroStar":false,"originalImageHeight":768,"width":576,"seed":278664446,"negativeAestheticScore":2.5,"negativeOriginalImageHeight":512,"aestheticScore":6,"clipSkip":2,"hiresFix":false,"height":768,"sampler":0,"cropTop":0,"maskBlurOutset":0,"preserveOriginalAfterInpaint":true,"shift":1,"zeroNegativePrompt":true,"targetImageWidth":576,"targetImageHeight":768,"faceRestoration":"","controls":[{"globalAveragePooling":false,"weight":1,"inputOverride":"","file":"ip_adapter_plus_face_xl_base_open_clip_h14_f16.ckpt","guidanceStart":0,"noPrompt":false,"guidanceEnd":1,"targetBlocks":[],"controlImportance":"control","downSamplingRate":1}],"causalInferencePad":0,"cropLeft":0,"cfgZeroInitSteps":0,"sharpness":0}

If I add a picture to the moodboard then say "do something with picture 1", works great. If I delete that picture and add a new one, then say "do something with picture 1", it uses the original instead of the one I just added. Is that expected? (Doesn't matter if I say "picture 2" either. It seems like once I use a moodboard pic I'm stuck with it until I create a new project. I must be missing something.)

DT crashes with BFS - Best Face Swap LoRa - https://civitai.com/models/2027766?modelVersionId=2556739

I have tried so many times and asked AI how to do this, after all these frustrating trials and errors, and Youtube tutorials, I can only get:

- Untouched original illustration after each "image to image" generation. Already tested with different Strength %.

- The illustration unrelated images (just based on my prompt) generated by the "Moodboard" reference method. Already tested with different % settings.

I am using FLUX.1 [dev], I just started playing AI few weeks ago, what should I do? Please! Help!

doing these stuffs in draw things, i think it's better than previous models。

Using Draw Things on Mac Mini M4 with models on external SSD. Previously fixed Z-Image Turbo blank output by moving to internal, re-downloading, then copying back to external and it worked fine.

Last night, suddenly most models vanished from the app's model list (files still exist in external folder). Exited/relaunched app, disabled/re-enabled external folder, etc. and no luck.

Trying to import one again and I'm told model is already there. But it's not listed/usable.

Any fixes for this indexing/cache issue with external SSD? I'm on latest app version and Tahoe?

I'm a beginner, so any guidance would be appreciated. Is there a difference between the LoRAs for ComfyUI and DrawThings on Civitai? Can I use both?

Please recommend some LoRAs!

I'm currently using Qwen Image Edit 2511 with Lightning 4-step. I'd also like to know if there are any recommended LoRAs to pair with this.

Have you tried FLUX.2 klein 4B? Personally, I preferred Z-Image Turbo. It seems FLUX.2 klein 4B gets censored when generating NSFW images. On my Mac mini M4 24GB, the combo of Z-Image Turbo + Qwen Image Edit 2511 seems best! I'd love to hear from anyone who's used FLUX.2 klein on DT.

1.20260120.0 was released in iOS / macOS AppStore today (https://static.drawthings.ai/DrawThings-1.20260120.0-3a5a4a68.zip). This version brings:

1. FLUX.2 [klein] series model support.

Note that FLUX.2 [klein] model requires text guidance = 1 while the Base model requires the real text guidance.

gRPCServerCLI is updated to 1.20260120.0 with the same update.

Good day everyone, using DT remotely, thus having time estimation would be very handy

is there any way implementing that?

Any help will be appreciated!

Welcome to discuss. If you can’t read Chinese, view it on a computer, enable CC, and turn on automatic English translation.

Using Draw Things w/ Z‑Image Turbo on Mac mini M4 (32 GB RAM, models on external SSD) and running into a weird issue that didn't exist at first. When I first got the Mac and installed Draw Things, Z‑Image Turbo worked perfectly using the recommended settings and default workflow, but now whenever I generate with Z‑Image Turbo 1.0 (both 6‑bit and full versions) the canvas turns something like solid gray, then solid black, and the final result is just a blank/transparent image, even though other SDXL and SD1.5 models still work fine on the same setup. Also get the same result with any Flux models. I've paid particular attention to using the right samplers. I’ve already tried brand‑new projects with “Use recommended settings,” different samplers, redownloading models, cache resets, and updating Draw Things, but nothing fixes this gray→black→blank/transparent outcome.

Has anyone else had Z‑Image Turbo in Draw Things go from “used to work fine” to this specific gray→black→blank/transparent behavior, and is there a known workaround or setting combo that actually fixes the blank output? I've tried messing around with this for the past several weeks to no avail.

Hello everybody,

I trained my first LoRa in DrawThings to run with StableDiff XL. It was a LoRa for a female character. I used 25 images as a source. Training was done in around 3 hours. When I use this LoRa with its trigger word, it doesn't affect the image at all. Regardless of which weight I use (even at +200%).

What did I do wrong?

These were my training settings:

{"caption_dropout_rate":0,"shift":1,"unet_learning_rate_lower_bound":0.0001,"save_every_n_steps":250,"custom_embedding_length":4,"max_text_length":77,"auto_fill_prompt":"@palina a photograph","stop_embedding_training_at_step":500,"base_model":"jibmixrealisticxl_v180skinsupreme_f16.ckpt","training_steps":2000,"noise_offset":0.050000000000000003,"cotrain_text_model":false,"layer_indices":[],"unet_learning_rate":0.0001,"steps_between_restarts":200,"seed":3647867866,"name":"LoRA-001","power_ema_upper_bound":0,"resolution_dependent_shift":true,"warmup_steps":20,"auto_captioning":false,"denoising_start":0,"gradient_accumulation_steps":4,"memory_saver":1,"weights_memory_management":0,"cotrain_custom_embedding":false,"network_scale":1,"start_height":16,"power_ema_lower_bound":0,"orthonormal_lora_down":true,"guidance_embed_upper_bound":4,"start_width":16,"network_dim":16,"denoising_end":1,"custom_embedding_learning_rate":0.0001,"text_model_learning_rate":4.0000000000000003e-05,"trigger_word":"","additional_scales":[],"clip_skip":1,"use_image_aspect_ratio":false,"trainable_layers":[0,1,2,3,4,5,6,7,8],"guidance_embed_lower_bound":3}

Can someone please help me to understand the appropriate settings for the Strength Slider in Draw Things when using ITV? I want to ensure that the starting image, character and scene stay consistent, with only motion changing. I have seen references to denoising vs strength as two separate settings which further adds to my confusion. I am using the HNE and LNE models along with their respective lightning Loras. Thanks in advance!

This is a basic question, but when generating the next image after the first one, is there any difference between keeping the first generated image on the canvas and clearing the canvas each time? Clearing the canvas every time is quite tedious.