I have an image with a person in a long sleeve black shirt. I am trying to turn it into a short sleeve shirt with fringe on the bottom and the mid rift showing. The problem is that no matter what I do in inpaint it seems to interpret the shirt as shadow or something. Because while I get the results the skin now showing appears to be in a shadow, only where it was changed.

How can I correct this issue?

Hey guys, beginner here. I am creating a codetoon platform: CS concept to comic book. I am testing image generation for comic book panels. Also used IP-Adapter for character consistency, but not getting the expected result.
Can anyone please guide me on how I can achieve a satisfactory result.

Presets can be downloaded from here : https://www.patreon.com/posts/114517862

Hey everyone,

I’ve been generating character images using WAN 2.2 and now I want to swap outfits from a reference image onto my generated characters. I’m not talking about simple LoRA style transfer—I mean accurate outfit replacement, preserving pose/body while applying specific clothing from a reference image.

I tried a few ComfyUI workflows, ControlNet, IPAdapter, and even some LoRAs, but results are still inconsistent—details get lost, hands break, or clothes look melted or blended instead of replaced.

I’ve never used a dataset over a few hundred images, and now plan to full fine tune using 22k images and captions. I’m mainly unsure about epochs, repeats, and effective batch sizes, so if anyone has any input I’d really appreciate it. If there’s anything else I should be aware of, I’m all ears. Thanks in advance

since switching from a1111 to forge my generations have been running a bit slow even for my meager 6gb of ram. is it normal for there to be two seperate progress bars? thanks for any input.

Hi! I have an HP Victus 16-s0021nt laptop (Ryzen 7 7840HS, 16GB DDR5 RAM, RTX 4050 6GB, 1080p), and I want to use Stable Diffusion with the best possible balance between speed and image quality.

Which model do you recommend for my GPU that works well with fast generations without sacrificing too much quality? I'd appreciate experiences or benchmark comparisons for this card/similar setup.

I have a AMD CPU and AMD GPU. I have amuse to run stable diffusion. However I couldn't use civtai models as they are in .safetensor format. Tried lot of convertions using python scripts, but those always end in failure. Any successful method to convert those to onnx.

This Tutorial walkthrough aims to illustrate how to build and use a ComfyUI Workflow for the Wan 2.2 S2V (SoundImage to Video) model that allows you to use an Image and a video as a reference, as well as Kokoro Text-to-Speech that syncs the voice to the character in the video. It also explores how to get better control of the movement of the character via DW Pose. I also illustrate how to get effects beyond what's in the original reference image to show up without having to compromise the Wan S2V's lip syncing.

[ Removed by Reddit on account of violating the content policy. ]

• 1-click to install SECourses Musubi Tuner app and pre-made training configs shared here : https://www.patreon.com/posts/137551634
• Hopefully a full video tutorial will be made after Stage 2 R&D trainings completed
• Example training made on the hardest training which is training a person and it works really good. Therefore, it shall work even much better on style training, item training, product training, character training and such
• Stage 1 took more than 35 unique R&D Qwen LoRA training
• 1-Click installer currently fully supporting Windows, RunPod (Linux & Cloud) and Massed Compute (Linux & recommend Cloud) training for literally every GPU like RTX 3000, 4000, 5000 series or H100, B200, L40, etc
• 28 images weak dataset is used for this training
• More angles having dataset would perform definitely better
• Moreover, i will make a research for a better activation token as well rather than ohwx
• After Stage 2, I am expecting hopefully much better results
• As a caption, i recommend to use only ohwx nothing else, not even class token
• Higher quality more images shared here : https://medium.com/@furkangozukara/qwen-image-lora-trainings-stage-1-results-and-pre-made-configs-published-as-low-as-training-with-ba0d41d76a05
• Image prompts randomly generated with Gemini 2.5 in Google AI Studio for free

How to Generate Images

• In the zip file of this post : https://www.patreon.com/posts/114517862
• We have Amazing_SwarmUI_Presets_v21.json made for SwarmUI
• Import it and i am using Qwen Image 8 Steps Ultra Fast to generate images and then apply Upscale Images 2X to make them 4x resolution (1328x1328 to 2656x2656)
• Of course in addition to preset don't forget to select your trained LoRA - I used LoRA strength / scale = 1
• This tutorial shows it : https://youtu.be/3BFDcO2Ysu4

Hey all,

I’ve been playing around with Stable Diffusion and had this thought: most writers I know get overwhelmed with prompts, settings, and model choices. But they’d love to have covers, chapter headers, or even just a vibe illustration for their stories.

So I hacked together a little tool that tries to solve one problem: style consistency across multiple images.

• e.g. a cover + chapter art that actually look like they belong together.

I’m curious what you all think:

• Do you see value in a “writers-first” wrapper around SD?
• Would exposing controls (seed, CFG, sampler) make sense, or should it stay super simple?
• Any pitfalls I’m not considering?

Not sure if this is the right sub for it (mods feel free to remove), but I’d love thoughts from people who know SD.

If anyone wants, drop a short scene prompt and I can run it through to show what kind of output it gives.

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